# Multi-Style Adaptation

Alma Photographic PDXL is a text-to-image generation model based on Stable Diffusion XL, specializing in photorealistic imagery with particular expertise in sci-fi, fantasy, monster, and pony themes.

A text-to-image generation model optimized based on Animagine XL 4.0, specifically enhanced for anime and cyberpunk styles.

A text-to-image model based on Stable Diffusion XL, specializing in anime-style image generation, particularly skilled at drawing aesthetically pleasing girl characters.

A text-to-image generation model based on Stable Diffusion XL, specializing in anime character and style creation, suitable for beginners.

A text-to-image generation model based on Stable Diffusion XL (SDXL), supporting ControlNet functionality.

A text-to-image generation model based on Stable Diffusion XL, supporting ControlNet functionality.

A text-to-image generation model based on Stable Diffusion XL (SDXL), supporting ControlNet functionality.

This is a LoRA model based on FLUX.1-dev and SDXL, specifically designed for generating images with Emma Watson's facial features.

EcomID is a text-to-image generation model based on SDXL, specializing in ID-customized generation through a single reference image while maintaining semantic consistency and facial keypoint control.

A depth control adapter based on the T2I-Adapter architecture, designed to enhance the depth perception capability of SDXL models in image generation

Flux-Dev2Pro is a text-to-image generation model based on the diffusers library, capable of producing high-quality images from textual descriptions.

A clothing-focused LoRA model based on the Florence-2 annotation system, specializing in generating high-quality fashion design images with support for multiple styles and material expressions

A collection of Stable Diffusion prompt templates for generating high-quality diverse images, covering surreal, architectural, photographic and other styles

A text-to-image generation model based on StableDiffusionXLPipeline, focused on artistic creation.

The Stable Diffusion Model is a text-to-image generation model capable of producing high-quality images based on textual descriptions.

ProteusV0.2 is an advanced version based on OpenDalleV1.1, achieving exceptional text-to-image generation through core feature enhancements, with significant improvements in prompt comprehension and style representation.

MAGNeT is a text-to-music and text-to-sound model capable of generating high-quality audio samples from text descriptions.

Animagine XL 3.0 is the latest open-source anime text-to-image model developed based on Stable Diffusion XL, focusing on high-quality anime image generation with improved hand anatomy and conceptual understanding.

The Stable Diffusion Model is a text-to-image diffusion model capable of generating high-quality images based on textual descriptions.

An art-style text-to-image model fine-tuned based on Stable Diffusion v1.5, excelling in blending anime and photorealistic styles

An art-style text-to-image model fine-tuned on Stable Diffusion XL 1.0, excelling in anime-style and artistic image generation

A text-to-image generation model based on diffusion models, capable of producing high-quality images from textual descriptions.

RealCartoon3D is a text-to-image generation model based on Stable Diffusion technology, specializing in generating images with a 3D cartoon style.

The Stable Diffusion Model is a text-to-image generation model capable of producing high-quality images based on textual descriptions.

The Stable Diffusion Model is a text-to-image generation model capable of producing high-quality images based on textual descriptions.

This is a text-to-image generation model based on stable diffusion technology, capable of producing high-quality images from text descriptions.

A text-to-image generation model based on stable diffusion technology, capable of producing high-quality images from textual descriptions.

MusicGen is a text-to-music model that generates high-quality music samples based on text descriptions or audio prompts.

An automatic prompt optimization model that generates professional-level prompts from minimal descriptions for high-quality image generation.

Anything V3 is a text-to-image generation model based on Stable Diffusion technology, capable of generating high-quality images from textual descriptions.

This is a fish image classification model built based on PyTorch and HuggingPics, capable of recognizing different types of 2D fish images.

A text-to-image generation model based on stable diffusion technology, focused on artistic creation

A Stable Diffusion model trained on anime characters from 'ONIMAI: I'm Now Your Sister!', capable of generating high-quality anime images of 6 main characters

A text-to-image model based on Stable Diffusion technology, capable of generating high-quality images from text descriptions

A text-to-image generation model based on diffusion models, capable of producing high-quality images from text descriptions

Wind v2 is a text-to-image generation model based on the Stable Diffusion architecture, released under the CreativeML OpenRAIL-M license.

Diffusers is a diffusion model library for image generation, supporting various pre-trained models and custom training.

Stable Diffusion is a diffusion model for text-to-image generation, capable of producing high-quality images based on textual descriptions.

A diffusion model for text-to-image generation, capable of producing high-quality images from textual descriptions.